模型:
runwayml/stable-diffusion-inpainting
Stable Diffusion Inpainting is a latent text-to-image diffusion model capable of generating photo-realistic images given any text input, with the extra capability of inpainting the pictures by using a mask.
The Stable-Diffusion-Inpainting was initialized with the weights of the Stable-Diffusion-v-1-2 . First 595k steps regular training, then 440k steps of inpainting training at resolution 512x512 on “laion-aesthetics v2 5+” and 10% dropping of the text-conditioning to improve classifier-free classifier-free guidance sampling . For inpainting, the UNet has 5 additional input channels (4 for the encoded masked-image and 1 for the mask itself) whose weights were zero-initialized after restoring the non-inpainting checkpoint. During training, we generate synthetic masks and in 25% mask everything.
You can use this both with the ?Diffusers library and the RunwayML GitHub repository .
from diffusers import StableDiffusionInpaintPipeline pipe = StableDiffusionInpaintPipeline.from_pretrained( "runwayml/stable-diffusion-inpainting", revision="fp16", torch_dtype=torch.float16, ) prompt = "Face of a yellow cat, high resolution, sitting on a park bench" #image and mask_image should be PIL images. #The mask structure is white for inpainting and black for keeping as is image = pipe(prompt=prompt, image=image, mask_image=mask_image).images[0] image.save("./yellow_cat_on_park_bench.png")
How it works:
image | mask_image |
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prompt | Output |
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Face of a yellow cat, high resolution, sitting on a park bench |
Developed by: Robin Rombach, Patrick Esser
Model type: Diffusion-based text-to-image generation model
Language(s): English
License: The CreativeML OpenRAIL M license is an Open RAIL M license , adapted from the work that BigScience and the RAIL Initiative are jointly carrying in the area of responsible AI licensing. See also the article about the BLOOM Open RAIL license on which our license is based.
Model Description: This is a model that can be used to generate and modify images based on text prompts. It is a Latent Diffusion Model that uses a fixed, pretrained text encoder ( CLIP ViT-L/14 ) as suggested in the Imagen paper .
Resources for more information: GitHub Repository , Paper .
Cite as:
@InProceedings{Rombach_2022_CVPR, author = {Rombach, Robin and Blattmann, Andreas and Lorenz, Dominik and Esser, Patrick and Ommer, Bj\"orn}, title = {High-Resolution Image Synthesis With Latent Diffusion Models}, booktitle = {Proceedings of the IEEE/CVF Conference on Computer Vision and Pattern Recognition (CVPR)}, month = {June}, year = {2022}, pages = {10684-10695} }
The model is intended for research purposes only. Possible research areas and tasks include
Excluded uses are described below.
Note: This section is taken from the DALLE-MINI model card , but applies in the same way to Stable Diffusion v1 .
The model should not be used to intentionally create or disseminate images that create hostile or alienating environments for people. This includes generating images that people would foreseeably find disturbing, distressing, or offensive; or content that propagates historical or current stereotypes.
Out-of-Scope UseThe model was not trained to be factual or true representations of people or events, and therefore using the model to generate such content is out-of-scope for the abilities of this model.
Misuse and Malicious UseUsing the model to generate content that is cruel to individuals is a misuse of this model. This includes, but is not limited to:
While the capabilities of image generation models are impressive, they can also reinforce or exacerbate social biases. Stable Diffusion v1 was trained on subsets of LAION-2B(en) , which consists of images that are primarily limited to English descriptions. Texts and images from communities and cultures that use other languages are likely to be insufficiently accounted for. This affects the overall output of the model, as white and western cultures are often set as the default. Further, the ability of the model to generate content with non-English prompts is significantly worse than with English-language prompts.
Training Data The model developers used the following dataset for training the model:
Training Procedure Stable Diffusion v1 is a latent diffusion model which combines an autoencoder with a diffusion model that is trained in the latent space of the autoencoder. During training,
We currently provide six checkpoints, sd-v1-1.ckpt , sd-v1-2.ckpt and sd-v1-3.ckpt , sd-v1-4.ckpt , sd-v1-5.ckpt and sd-v1-5-inpainting.ckpt which were trained as follows,
sd-v1-1.ckpt : 237k steps at resolution 256x256 on laion2B-en . 194k steps at resolution 512x512 on laion-high-resolution (170M examples from LAION-5B with resolution >= 1024x1024 ).
sd-v1-2.ckpt : Resumed from sd-v1-1.ckpt . 515k steps at resolution 512x512 on "laion-improved-aesthetics" (a subset of laion2B-en, filtered to images with an original size >= 512x512 , estimated aesthetics score > 5.0 , and an estimated watermark probability < 0.5 . The watermark estimate is from the LAION-5B metadata, the aesthetics score is estimated using an improved aesthetics estimator ).
sd-v1-3.ckpt : Resumed from sd-v1-2.ckpt . 195k steps at resolution 512x512 on "laion-improved-aesthetics" and 10% dropping of the text-conditioning to improve classifier-free guidance sampling .
sd-v1-4.ckpt : Resumed from stable-diffusion-v1-2.225,000 steps at resolution 512x512 on "laion-aesthetics v2 5+" and 10 % dropping of the text-conditioning to classifier-free guidance sampling .
sd-v1-5.ckpt : Resumed from sd-v1-2.ckpt. 595k steps at resolution 512x512 on "laion-aesthetics v2 5+" and 10% dropping of the text-conditioning to improve classifier-free guidance sampling.
sd-v1-5-inpaint.ckpt : Resumed from sd-v1-2.ckpt. 595k steps at resolution 512x512 on "laion-aesthetics v2 5+" and 10% dropping of the text-conditioning to improve classifier-free guidance sampling. Then 440k steps of inpainting training at resolution 512x512 on “laion-aesthetics v2 5+” and 10% dropping of the text-conditioning. For inpainting, the UNet has 5 additional input channels (4 for the encoded masked-image and 1 for the mask itself) whose weights were zero-initialized after restoring the non-inpainting checkpoint. During training, we generate synthetic masks and in 25% mask everything.
Hardware: 32 x 8 x A100 GPUs
Optimizer: AdamW
Gradient Accumulations : 2
Batch: 32 x 8 x 2 x 4 = 2048
Learning rate: warmup to 0.0001 for 10,000 steps and then kept constant
Evaluations with different classifier-free guidance scales (1.5, 2.0, 3.0, 4.0, 5.0, 6.0, 7.0, 8.0) and 50 PLMS sampling steps show the relative improvements of the checkpoints:
Evaluated using 50 PLMS steps and 10000 random prompts from the COCO2017 validation set, evaluated at 512x512 resolution. Not optimized for FID scores.
To assess the performance of the inpainting model, we used the same evaluation protocol as in our LDM paper . Since the Stable Diffusion Inpainting Model acccepts a text input, we simply used a fixed prompt of photograph of a beautiful empty scene, highest quality settings .
Model | FID | LPIPS |
---|---|---|
Stable Diffusion Inpainting | 1.00 | 0.141 (+- 0.082) |
Latent Diffusion Inpainting | 1.50 | 0.137 (+- 0.080) |
CoModGAN | 1.82 | 0.15 |
LaMa | 2.21 | 0.134 (+- 0.080) |
Stable Diffusion v1 Estimated Emissions Based on that information, we estimate the following CO2 emissions using the Machine Learning Impact calculator presented in Lacoste et al. (2019) . The hardware, runtime, cloud provider, and compute region were utilized to estimate the carbon impact.
@InProceedings{Rombach_2022_CVPR, author = {Rombach, Robin and Blattmann, Andreas and Lorenz, Dominik and Esser, Patrick and Ommer, Bj\"orn}, title = {High-Resolution Image Synthesis With Latent Diffusion Models}, booktitle = {Proceedings of the IEEE/CVF Conference on Computer Vision and Pattern Recognition (CVPR)}, month = {June}, year = {2022}, pages = {10684-10695} }
This model card was written by: Robin Rombach and Patrick Esser and is based on the DALL-E Mini model card .