模型:
CompVis/ldm-celebahq-256
Paper : High-Resolution Image Synthesis with Latent Diffusion Models
Abstract :
By decomposing the image formation process into a sequential application of denoising autoencoders, diffusion models (DMs) achieve state-of-the-art synthesis results on image data and beyond. Additionally, their formulation allows for a guiding mechanism to control the image generation process without retraining. However, since these models typically operate directly in pixel space, optimization of powerful DMs often consumes hundreds of GPU days and inference is expensive due to sequential evaluations. To enable DM training on limited computational resources while retaining their quality and flexibility, we apply them in the latent space of powerful pretrained autoencoders. In contrast to previous work, training diffusion models on such a representation allows for the first time to reach a near-optimal point between complexity reduction and detail preservation, greatly boosting visual fidelity. By introducing cross-attention layers into the model architecture, we turn diffusion models into powerful and flexible generators for general conditioning inputs such as text or bounding boxes and high-resolution synthesis becomes possible in a convolutional manner. Our latent diffusion models (LDMs) achieve a new state of the art for image inpainting and highly competitive performance on various tasks, including unconditional image generation, semantic scene synthesis, and super-resolution, while significantly reducing computational requirements compared to pixel-based DMs.
Authors
Robin Rombach, Andreas Blattmann, Dominik Lorenz, Patrick Esser, Björn Ommer
!pip install diffusers from diffusers import DiffusionPipeline model_id = "CompVis/ldm-celebahq-256" # load model and scheduler pipeline = DiffusionPipeline.from_pretrained(model_id) # run pipeline in inference (sample random noise and denoise) image = pipeline(num_inference_steps=200)["sample"] # save image image[0].save("ldm_generated_image.png")
!pip install diffusers from diffusers import UNet2DModel, DDIMScheduler, VQModel import torch import PIL.Image import numpy as np import tqdm seed = 3 # load all models unet = UNet2DModel.from_pretrained("CompVis/ldm-celebahq-256", subfolder="unet") vqvae = VQModel.from_pretrained("CompVis/ldm-celebahq-256", subfolder="vqvae") scheduler = DDIMScheduler.from_config("CompVis/ldm-celebahq-256", subfolder="scheduler") # set to cuda torch_device = "cuda" if torch.cuda.is_available() else "cpu" unet.to(torch_device) vqvae.to(torch_device) # generate gaussian noise to be decoded generator = torch.manual_seed(seed) noise = torch.randn( (1, unet.in_channels, unet.sample_size, unet.sample_size), generator=generator, ).to(torch_device) # set inference steps for DDIM scheduler.set_timesteps(num_inference_steps=200) image = noise for t in tqdm.tqdm(scheduler.timesteps): # predict noise residual of previous image with torch.no_grad(): residual = unet(image, t)["sample"] # compute previous image x_t according to DDIM formula prev_image = scheduler.step(residual, t, image, eta=0.0)["prev_sample"] # x_t-1 -> x_t image = prev_image # decode image with vae with torch.no_grad(): image = vqvae.decode(image) # process image image_processed = image.cpu().permute(0, 2, 3, 1) image_processed = (image_processed + 1.0) * 127.5 image_processed = image_processed.clamp(0, 255).numpy().astype(np.uint8) image_pil = PIL.Image.fromarray(image_processed[0]) image_pil.save(f"generated_image_{seed}.png")